R stable diffusion.

Stable Diffusion tagging test. This is the Stable Diffusion 1.5 tagging matrix it has over 75 tags tested with more than 4 prompts with 7 CFG scale, 20 steps, and K Euler A sampler. With this data, I will try to decrypt what each tag does to your final result. So let's start: Discussion. Curious to know if everyone uses the latest Stable Diffusion XL engine now or if there are pros and cons to still using older engines vs newer ones. When using the API, what other API do you tend to use all the available parameters to optimise image generation or just stick with prompts, steps and width/height? Sort by: Add a ... This was very useful, thanks a lot for posting it! I was mainly interested in the painting Upscaler, so I conducted a few tests, including with two Upscalers that have not been tested (and one of them seems better than ESRGAN_4x and General-WDN. 4x_foolhardy_Remacri with 0 denoise, as to perfectly replicate a photo. I use MidJourney often to create images and then using the Auto Stable Diffusion web plugin, edit the faces and details to enhance images. In MJ I used the prompt: movie poster of three people standing in front of gundam style mecha bright background motion blur dynamic lines --ar 2:3

Generating iPhone-style photos. Most pictures I make with Realistic Vision or Stable Diffusion have a studio lighting feel to them and look like professional photography. The person in the foreground is always in focus against a blurry background. I’d really like to make regular, iPhone-style photos. Without the focus and studio lighting.

I created a reference page by using the prompt "a rabbit, by [artist]" with over 500+ artist names. It serves as a quick reference as to what the artist's style yields. Notice there are cases where the output is barely recognizable as a rabbit. Others are delightfully strange. It includes every name I could find in prompt guides, lists of ... It's late and I'm on my phone so I'll try to check your link in the morning. One thing that really bugs me is that I used to live "X/Y" graph because it I set the batch to 2, 3, 4 etc images it would show ALL of them on the grid png not just the first one. I assume there must be a way w this X,Y,Z version, but everytime I try to have it com

It won't let you use multiple GPUs to work on a single image, but it will let you manage all 4 GPUs to simultaneously create images from a queue of prompts (which the tool will also help you create). Just made the git repo public today after a few weeks of testing. There are probably still some issues but I've been running it on a 3 GPU rig 24/ ...Stable Diffusion 3 combines a diffusion transformer architecture and flow matching. We will publish a detailed technical report soon. We believe in safe, … This sometimes produces unattractive hair styles if the model is inflexible. But for the purposes of producing a face model for inpainting, this can be acceptable. HardenMuhPants. • 10 mo. ago. Just to add a few more simple terms style hair cuts. Whispy updo. Although these images are quite small, the upscalers built into most versions of Stable Diffusion seem to do a good job of making your pictures bigger with options to smooth out flaws like wonky faces (use the GFPGAN or codeformer settings). This is found under the "extras" tab in Automatic1111 Hope that makes sense (and answers your question). This was very useful, thanks a lot for posting it! I was mainly interested in the painting Upscaler, so I conducted a few tests, including with two Upscalers that have not been …

Osmosis is an example of simple diffusion. Simple diffusion is the process by which a solution or gas moves from high particle concentration areas to low particle concentration are...

Stable diffusion vs Midjourney. You can do it in SD aswell, but it requires far more efdort. Basically a lot of inpainting. Use custom models OP. Dream like and open journey are good ones if you like midjourny style. You can even train your own custom model with whatever style you desire. As I have said stable is a god at learning.

Stable Diffusion v2 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text …I'm able to get pretty good variations of photorealistic people using "contact sheet" or "comp card" in my prompts. But I'm also trying to use img2img to get a consistent set of different crops, expressions, clothing, backgrounds, etc, so any model or embedding I ...Someone told me the good images from stable diffusion are cherry picked one out hundreds, and that image was later inpainted and outpainted and refined and photoshoped etc. If this is the case the stable diffusion if not there yet. Paid AI is already delivering amazing results with no effort. I use midjourney and I am satisfied, I just wante ... Go to your Stablediffusion folder. Delete the "VENV" folder. Start "webui-user.bat". it will re-install the VENV folder (this will take a few minutes) WebUI will crash. Close Webui. now go to VENV folder > scripts. click folder path at the top. type CMD to open command window. I have been long curious about the popularity of Stable Diffusion WebUI extensions. There are so many extensions in the official index, many of them I haven't explore. Today, on 2023.05.23: I gathered the Github stars of all extensions in the official index.Negatives: “in focus, professional, studio”. Do not use traditional negatives or positives for better quality. MuseratoPC. •. I found that the use of negative embeddings like easynegative tends to “modelize” people a lot, makes them all supermodel photoshop type images. Did you also try, shot on iPhone in your prompt?

Stable Diffusion is a latent text-to-image diffusion model. Thanks to a generous compute donation from Stability AI and support from LAION, we were able to train a Latent … Steps for getting better images. Prompt Included. 1. Craft your prompt. The two keys to getting what you want out of Stable Diffusion are to find the right seed, and to find the right prompt. Getting a single sample and using a lackluster prompt will almost always result in a terrible result, even with a lot of steps. Hello everyone, I'm sure many of us are already using IP Adapter. But recently Matteo, the author of the extension himself (Shoutout to Matteo for his amazing work) made a video about character control of their face and clothing. Realistic. 38 10. u/Negative-Use274. • 16 hr. ago. NSFW. Something Sweet. Realistic. 28 2. r/sdnsfw: This sub is for all those who want to enjoy the new freedom that AI offers us to the fullest and without censorship. If you want to try Stable Diffusion v2 prompts, you can have a free account here (don't forget to choose SD 2 engine) https://app.usp.ai. The prompt book is showing different examples based on the official guide, with some tweaks and changes. Since it is using multi prompting and weights, use it for Stable Diffusion 2.1 up.

What are currently the best stable diffusion models? "Best" is difficult to apply to any single model. It really depends on what fits the project, and there are many good choices. CivitAI is definitely a good place to browse with lots of example images and prompts. I keep older versions of the same models because I can't decide which one is ...

This sometimes produces unattractive hair styles if the model is inflexible. But for the purposes of producing a face model for inpainting, this can be acceptable. HardenMuhPants. • 10 mo. ago. Just to add a few more simple terms style hair cuts. Whispy updo. Text-to-image generation is still on the works, because Stable-Diffusion was not trained on these dimensions, so it suffer from coherence. Note: In the past, generating large images with SD was possible, but the key improvement lies in the fact that we can now achieve speeds that are 3 to 4 times faster, especially at 4K resolution. This shift ...Fixing excessive contrast/saturation resulting from high CFG scales. At high CFG scales (especially >20, but often below that as well), generated images tend to have excessive and undesired levels of contrast and saturation. This is worse when using certain samplers and better when using others (from personal experience, k_euler is the best ...One of the most criticized aspects of cryptocurrencies is the fact that they change in value dramatically over short periods of time. Imagine you bought $100 worth of an ICO’s toke... Discussion. Curious to know if everyone uses the latest Stable Diffusion XL engine now or if there are pros and cons to still using older engines vs newer ones. When using the API, what other API do you tend to use all the available parameters to optimise image generation or just stick with prompts, steps and width/height? Sort by: Add a ... Skin color options were determined by the terms used in the Fitzpatrick Scale that groups tones into 6 major types based on the density of epidermal melanin and the risk of skin cancer. The prompt used was: photo, woman, portrait, standing, young, age 30, VARIABLE skin. Skin Color Variation Examples.installing stable diffusion. Hi, everyone, I have tried for weeks to figure out a way to download and run stable diffusion, but I can't seem to figure it out. Could someone point … Discussion. Curious to know if everyone uses the latest Stable Diffusion XL engine now or if there are pros and cons to still using older engines vs newer ones. When using the API, what other API do you tend to use all the available parameters to optimise image generation or just stick with prompts, steps and width/height? Sort by: Add a ... By default it's looking in your models folder. I needed it to look one folder deeper to stable-diffusion-webui\models\ControlNet. I think some tutorials are also having you put them in the stable-diffusion-webui\extensions\sd-webui-controlenet>models folder. Copy path and paste 'em in wherever you're saving 'em.

Hi. Below, I present my results using this tutorial. The original image (512x768) was created in Stable Diffusion (A1111), transferred to Photopea, resized to 1024x1024 (white background), and retransferred to txt2img (with original image prompt) using ControlNet ...

You can use agent scheduler to avoid having to use disjunctives by queuing different prompts in a row. Prompt S/R is one of more difficult to understand modes of operation for X/Y Plot. S/R stands for search/replace, and that's what it does - you input a list of words or phrases, it takes the first from the list and treats it as keyword, and ...

Negatives: “in focus, professional, studio”. Do not use traditional negatives or positives for better quality. MuseratoPC. •. I found that the use of negative embeddings like easynegative tends to “modelize” people a lot, makes them all supermodel photoshop type images. Did you also try, shot on iPhone in your prompt?In stable diffusion Automatic1111. Go to Settings tab. On the left choose User Interface. Then search for [info] Quicksettings list, by default you already should have sd_model_checkpoint there in the list, so there, you can add the word tiling. Go up an click Apply Settings then on Reload UI. After reload on top next to checkpoint you should ...If you want to try Stable Diffusion v2 prompts, you can have a free account here (don't forget to choose SD 2 engine) https://app.usp.ai. The prompt book is showing different examples based on the official guide, with some tweaks and changes. Since it is using multi prompting and weights, use it for Stable Diffusion 2.1 up.Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and then finetuned on 512x512 images. Note: Stable Diffusion v1 is a general text-to-image diffusion ... In the context of Stable Diffusion, converging means that the model is gradually approaching a stable state. This means that the model is no longer changing significantly, and the generated images are becoming more realistic. There are a few different ways to measure convergence in Stable Diffusion. So it turns out you can use img2img in order to make people in photos look younger or older. Essentially add "XX year old man/women/whatever", and set prompt strength to something low (in order to stay close to the source). It's a bit hit or miss and you probably want to run face correction afterwards, but itworks. 489K subscribers in the ... Stable Diffusion is a deep learning model used for converting text to images. It can generate high-quality, photo-realistic images that look like real photographs by simply inputting any text. The latest version of this model is Stable Diffusion XL, which has a larger UNet backbone network and can generate even higher quality images. Stable Diffusion XL Benchmarks. A set of benchmarks targeting different stable diffusion implementations to have a better understanding of their performance and scalability. Not surprisingly TensorRT is the fastest way to run Stable Diffusion XL right now. Interesting to follow if compiled torch will catch up with TensorRT.Use one or both in combination. The more information surrounding the face that SD has to take into account and generate, the more details and hence confusion can end up in the output. With focus on the face that’s all SD has to consider, and the chance of clarity goes up. bmemac. • 2 yr. ago.

Tesla M40 24GB - half - 31.64s. Tesla M40 24GB - single - 31.11s. If I limit power to 85% it reduces heat a ton and the numbers become: NVIDIA GeForce RTX 3060 12GB - half - 11.56s. NVIDIA GeForce RTX 3060 12GB - single - 18.97s. Tesla M40 24GB - half - 32.5s. Tesla M40 24GB - single - 32.39s. Stable Diffusion for AMD GPUs on Windows using DirectML. SD Image Generator. Simple and easy to use program. Lama Cleaner. One click installer in-painting tool to ... Hello, Im a 3d charactrer artist, and recently started learning stable diffusion. I find it very useful and fun to work with. Im still a beginner, so I would like to start getting into it a bit more. r/StableDiffusion. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. More info: …Instagram:https://instagram. pollen count patchoguebosphorus straps848 667 2845spectrum news 9 albany ny Code from Himuro-Majika's Stable Diffusion image metadata viewer browser extension \r"," Reading metadata with ExifReader, extra search results supported by String-Similarity \r"," Lazyload Script from Verlok, webfont is Google's Roboto, SVG icons from reds game final scorestar wars the clone wars revan fanfiction Someone told me the good images from stable diffusion are cherry picked one out hundreds, and that image was later inpainted and outpainted and refined and photoshoped etc. If this is the case the stable diffusion if not there yet. Paid AI is already delivering amazing results with no effort. I use midjourney and I am satisfied, I just wante ...This is an answer that someone corrects. The the base model seem to be tuned to start from nothing, then to get an image. The refiner refines the image making an existing image better. You can use the base model by it's self but for additional detail you should move to the second. Here for the answer. the super mario bros. movie showtimes near flint west 14 Making Stable Diffusion Results more like Midjourney. I was introduced to the world of AI art after finding a random video on YouTube and I've been hooked ever since. I love the images it generates but I don't like having to do it through Discord and the limitation of 25 images or having to pay. So I did some research looking for AI Art that ...Description. Artificial Intelligence (AI)-based image generation techniques are revolutionizing various fields, and this package brings those capabilities into the R environment. In other words, it's not quite multimodal (Finetuned Diffusion kinda is, though. Wish there was an updated version of it). The basic demos online on Huggingface don't talk to each other, so I feel like I'm very behind compared to a lot of people.